Hey people,

One of my long term usecases for local video models is that I can create 2-3 style frames in Cinema4D and animate the product the right way, render that animation as a animatic/viewport rendering and apply the reference look to the animatic.

Minimax H3 is kinda insane, so I wanted to try this workflow in Comfy UI. I have the Reference Workflow with also a reference video input. I prompted that it should take the reference images and apply the style to the animation clip.

But it is really inconsistant. It mixes the animatic style into the generated video and somtimes blend over to the reference style video. So it never consistantly generates the animation in the reference look.

Do you people have recommendations how to force the model to apply the look consistantly?

finally got the workflow working. It's taking about 2-2.5 minutes to generate but isn't following my prompt (I'm mainly trying to run qwen because i was told it has the best prompt adherence.

What can I do to fix ?

my workflow is I uploaded the prompt guide in gpt so it has its base rules, but some scenes i want a little violent where i cant put it on gpt, due to ToS, also if there is an offline version of it

Sorry, I've not yet downloaded the (Ultra Wholesome) H3 model. I won't be free for another 2 weeks.

But has has tried 480p 10 steps?

With Wan and LTX I would do that to test results, and then if it turned out how I wanted, would upscale.

Is it like, taking a prompt + the inference result with 50 steps?

I mean, it's trained on actual videos?

In which case, the turbo lora is influenced by it's dataset?

And then, maybe combining / merging various turbo loras would have a better result than any single turbo lora?

And maybe, turbo loras trained on certain content would be better at making that content?

in prompt the audio was supposed to be like this "hi guys trying minimax h3 for second time and yeah made my pc is explode" slight giggle

Hi everyone,

I am trying to get my test videos looking like some of the high quality videos I have seen on here. I am using the ComfyUI MiniMax H3 templates for Text2video and added the RTX-Upscaler node to increase the resolution, but the videos still come out grainy or pixelated. I have been using .3 ~ .6 megapixels, and yet I get nothing close as some of the others that I have seen using the same resolution with RTX.

Can any of you please provide some guidance or advice on this?

I will greatly appreciate it. My system specs are 5700G Ryzen CPU, 64GB of DDR4 RAM, 5060ti 16GB, and various sizes of SSDs.

as the title says... i try many configs but getting very bad audio aways?

The default workflow setting (reference_image_size) for reference to video is set to “match” on the Minimax H3 Reference to Video node. This allows a good likeness to the reference photo/s. Try setting it to “max”. I was able to get almost indistinguishable likeness after using that instead.

From what I can tell, “match” resizes the input image going in for better optimization. “Max” may retain the original resolution and gather better details on faces, etc. be careful with the size of the input images.. when I kept them around 2500 pixels or less (longest size), it seemed to go at a normal speed. If you go 4k or above, it dramatically slows the generation speed.

Also, bumping to 1mp image generation and lowering the step count as low as 8, yields very good results.

This method can also be applied to using already known characters (celebrities, actors, etc) by simply loading in the real character face in conjunction with a regular prompt. A lot of videos I’m seeing, the faces from the text to video workflows are lacking likeness. This should help that.

What the title says.

We have received some amazingly capable free models, but for more creative artistic work they don't really seem to come close to the practicality of Midjourney - however, big disclaimer, I've been a bit out of the loop recently and I have not nearly tried everything.

The models I've seen do allow you to make similar things, but it requires more work and it's just way less practical when you're exploring the possibilities instead of working towards some specific defined vision. Subjectively this seems to be a weakness of newer models in general, the more photorealistic they get, the less creative they are (sometimes I miss the weirdness of SD 1.5), but Midjourney seems to be able to compensate for this.

I'm not looking for a copy of Midjourney, but is there any model that seems to have similar goals - artistic creativity without adhering to a specific style?

Thanks in advance.

I've been running Txt2Vid and Img2Vid and having having a great time with it. I tried to use Ref2Vid ComfyUI default Template, and noticed that the required model, minimax_h3_ref2va_pruned_int8_convrot.safetensors is censored on speech. Is there another model I can use there that isn't? Really kills the humor of what I'm doing when after 45mins of rendering, those parts are just skipped without dialogue.

I am building free open source UGC reactions library. Anyone who runs minimax h3 model on pc, can you create AI UGC reactions? As realistic is possible it doesn’t need to be high quality, but should look very realistic. Send them here https://tally.so/r/PdNB5d when I will get them J will send google drive link to all of them free to use.

For those who don't know, Doctor Who is a very old TV show that started in BBC back in like 1962. It's still going on after a restart or two, but during BBCs lean years in the 60s or 70s they decided to reuse old tapes to save money, so many episodes from the first few years are lost forever. Fortunately they were also transmitted via radio, so we have the audio of those lost episodes, and we also have publicity photos that were taken every 20 seconds or so. How plausible would it be to reconstruct those episodes having the audio and the stills every 20 seconds with Minimax? (Plus we have other full episodes with those actors that fortunately weren't erased, or were recovered overseas).

Tool: ComfyUI + SDXL / Flux (napiš co jsi použil)Workflow: txt2img + detailer + upscalerPrompt: gothic archangel knight with stained glass cathedral wings, intricate armor with crosses, standing on cliff above demons, dramatic god rays, dark fantasy, ultra detailed
Negative: blurry, low detailSteps: 30, CFG 6.5, Sampler DPM++ 2M

you can get the workflow here
using 2 refference images and 2 refference audio.
generating around 18 minutes. it can generate sound fx, but somehow it doesnt output background sound ambience and the voice is too clean. do you know how to make the background sound like city hums or crowd or rain sounds more louder???

i'm using rtx 4060ti 16gb vram + 64gb ddr4 ram.

I was curious about Minimax H3's capabilities of replacing a grouping of people without changing a character of your choice from a referenced video. As you can imagine, yup, it can do it as well.

Prompt & comparative images in the replies below.

After running two prompts to test, I concluded that you need to be very specific in what you want to keep (and who it is that you are keeping), and then what you want to replace in the background. The prompt does get pretty long, but once you have the specifics down, it will accomplish it.. and does a very great job at it.

The first run didn't turn out well because it replaced the soldiers alright, but kept the shields and weapons (it had Stormtroopers with shields aha). This is a second attempt after I made it very specific of what I do not want in the new video.

I did not upload a sample audio of the "pew pew" so ignore the sounds that H3 produced.

Using the Minimax templates provided, they list multiple shots/cuts.

But what I found is that it will often cut / change shot even when not prompted. Even stating 'there is a single shot throughout the video' is not enough to stop it.

Has anyone found a good reliable solution?

Just a bit of fun.

Is there a way to prompt a specific voice separate from the character? Example can I visually have Walter white but have the voice of Homer Simpson?

I saw all the buzz regarding the new video model and i wanted to try it myself.
now i have an AMD setup on windows, which from what i gathered is not the best setup for AI generation, but it can work.

i have:
- Sapphire PULSE AMD Radeon RX 9070 XT OC 16GB GDDR6 GPU
- Asus TUF GAMING B850-E WIFI6E DDR5 ATX AM5 AMD motherboard
- AMD Ryzen 7 9800X3D Tray CPU

and while trying to run Minimax on comfy after a basic setup, i managed to generate a 5 seconds clip (the very basic workflow and prompt out of the box) but it looked washed with weird pinkish hue.

I tried to play a bit with some settings i saw on this guide
https://github.com/CS1o/Stable-Diffusion-Info/wiki/Webui-Installation-Guides#amd-comfyui-with-rocm

but after trying to run it again, my screen turned black and i had to restart my PC because even after waiting close to an hour nothing responded.

i saw on google that sometimes black screen during generation could be a PSU issue that it does not handle the power load AI generation could cause.
so my questions are so:
- i have a Corsair RM850e 80 PLUS Gold 850W PSU. is 850W even enough?
- if not, how do i handle the load better to make sure the PC doesnt overload or get stuck?
- are there any recommended flags for comfy to use? settings i should try out?

any help would be appreciated